Why do we see rainbow colours on a CD?

Display of colours by a CD can be understood in terms of the working of a plane optical reflection grating. It is a flat optical device whose surface is ruled (striped) with a set of closely and uniformly spaced lines, such that light is reflected by the gaps and absorbed by the lines. , When light falls on a plane reflection grating, it is scattered in all directions by each of its reflecting stripes. These waves from individual gaps are termed wavelets. When we look at the grating from a distance, wavelets from different stripes travel different distances to reach the retina of our eye. Their crests (or troughs) do not reach a given point at the same time. Usually crests of some wavelets and troughs of others reach a point. Troughs have the property of partially or totally nullifying crests and vice versa, depending on their strengths. In this case the wavelets are said to interfere destructively.

But for certain orientations of the grating, it so happens that troughs (or crests) of all wavelets reach a point together, enhancing the effect of each other. The wavelets are then said to interfere constructively. Ordinarily, when light falls on a grating, wavelets corresponding to all wavelengths (or colours) are sent out by the reflecting gaps. Since the conditions of constructive interference hold good only for some particular wavelength, light intensity at the receiving point is exceptionally high only for that wavelength. Light from the grating from the related direction is thus rich in the corresponding colour. Similarly light from slightly different direction is rich in another colour. A CD has a data recording track, which spirals from its outer periphery to the inner circular boundary. This takes many tens of thousand rounds about the CD's centre. When examined along a radius of the CD, it is found to have a structure similar to that of a reflection grating — a set of almost straight tracks running perpendicular to the radius and separated by gaps. Therefore, like a grating the CD also displays colours. Courtesy : The Hindu

Why are diesel driven vehicles noisy?

Diesel engines are inherently noisy because of the auto-ignition of the initially formed mixture of fuel-vapour and air, which causes rapid rate of pressure rise producing the characteristic noise. This abnormality is due to the inborn feature of a diesel engine, which uses a high compression ratio to obtain high fuel efficiency and a high compression-temperature for ignition of the fuel injected into the cylinder at high pressure. In a diesel engine, unlike a petrol engine, the air and fuel do not mix outside the engine.

Air is compressed to a high compression ratio leading to high fuel efficiency, and also to a high compression temperature to initiate combustion. They just meet inside the combustion chamber for a brief period, while all the processes of mixture formation like fuel-jet break-up, evaporation and mixing should take place within a short time called ignition delay period. Combustion follows this after the initially formed fuel-air vapour auto-ignites with a noise, forming the sources of ignition (chemical spark plugs if you wish to call) for the bulk of the fuel remaining. Diesel engine can be made less noisy by using common rail high-pressure injection system and electronic control. Courtesy : The Hindu

Why does a mushroom shaped cloud form after a nuclear bomb explosion?

ANSWER I: When a nuclear weapon explodes, there is a rapid release of a large amount of energy within a small volume. This results in significant increase in temperature and pressure. The temperature may be a few tens of million degrees and pressure a few million times the atmospheric pressure. At this temperature, all the material present in the weapon will be converted into hot compressed gases.

Within a fraction of a millionth of a second of the explosion, the weapon's residues emit large amounts of energy mainly in the form of X-rays. The surrounding atmosphere absorbs this energy. This results in the formation of a blazing, highly luminous, spherical mass of air and gaseous weapon residues called the fireball.

Within an extremely short time after the explosion, the fire ball from a high yield nuclear weapon will be about 130 metres across increasing to about 1700 metres in ten seconds.

The fireball expands rapidly engulfing the surrounding air. The ball of hot air is less dense than the surrounding air. It rises swiftly like a hot air balloon.

This rising column pulls up debris of the weapon, dust and moisture along with it forming a cloud. As it moves up, it cools gradually and reaches about 10 km where the atmosphere is extremely stable.

The ball of air mass moving up does not have enough energy to penetrate this stable layer. It flattens out. As the relatively warmer layers at the bottom push up, the top layers spread laterally and equally in all directions, and the cooler denser layers descend at the edges, giving a distinct mushroom shape.

ANSWER II: Atmospheric nuclear explosion leads to sudden formation of a massive fireball near the ground, setting aflame whatever is in its vicinity. Since the fireball is very hot and thus less dense than the surrounding air, it rises very quickly. The massive updraft due to the rapidly rising fireball leaves a column of low-pressure. This acts as a chimney, sucking in smoke, dust and debris from the surroundings. This forms the stem of the mushroom.

At first the mixture of hot air and dust rises vertically, forming the column of the cloud. But as the hot cloud meets the colder air at higher altitudes, it slowly cools. Eventually the cloud reaches the temperature of the surrounding air and ceases to rise, but spreads horizontally along air levels at the same altitude, which are at the same temperature. This forms the cap of the mushroom.

The smoke, dust and debris gushing into the central column cause toroidal eddy currents in the horizontally spreading hot cloud. This introduces curling under the cap of the mushroom. Mushroom clouds are most commonly associated with nuclear weapons. However, any massive explosion capable of creating the same conditions would produce a mushroom cloud. Volcanic eruptions are typical natural mushroom clouds. Courtesy : The Hindu

Why are we unable to walk straight and tend to lose our balance when we walk with our eyes closed?

ANSWER I: When one wants to walk in a straight line the brain takes a visual reference point and uses this reference point to maintain the movement in a straight line. In the absence of a visual reference point the brain is unable to ascertain the exact direction in which the body is moving and hence most of us will be unable to walk in a straight line with our eyes closed.

However visual reference is not the only mechanism of maintaining balance or direction in our movements. We can judge the direction of movement very accurately with the balance sense from the balance mechanism in our inner ear and by kinesthetic sense, which is sensations we derive from our muscles, joints and ligaments.

Both these sensations will tell us very accurately as to the balance and direction of movement and position of the body at any point of time. This is highly developed in acrobats, gymnasts and ice skaters. They would be able walk in a straight line even with their eyes closed.

ANSWER II: Human balance is maintained by three pillars namely vestibular system — balance organ in the inner ear and its connections — eye and proprioceptors in the joints of the body. Man can maintain balance with any two pillars but not with one pillar alone. Imagine a blind man who can walk without losing balance.

If the same person develops any disorder in the vestibular system or proprioceptive system, then he cannot stand or walk. Similarly if a normal person develops a disorder in the inner ear or proprioceptive system and tries to stand with eyes closed and feet together, he will lose balance. Courtesy : The Hindu

Why do pregnant women like eating tamarind and raw mangoes?

Craving for unusual food (and some non food items like ash called pica) is considered as the first sign of pregnancy. In reality, although some women do get strong cravings, many do notNo one knows for sure what causes food cravings. Many women find that their senses of taste and smell are changed by pregnancy. For example, some women experience an odd metallic taste in their mouths very early in pregnancy (maybe the first sign of pregnancy for an `experienced' mother); others find that taste and smell are dulled. It is possible that these changes affect food likes and dislikes.

Some people think that cravings happen in response to temporary deficiency of specific nutrients. There is probably some truth in this, but it is not the whole story. We only need minute quantities of each vitamin and mineral — certainly not enough to justify a continual craving for just one food. For some women, food cravings may be a conscious or subconscious response to emotion.

They may crave a favourite childhood food, or a food that is of special significance to their religion or culture. Craving unusual foods may also be a private way of marking the special state of being pregnant.

Rather than develop a food craving, many women find they suddenly go off certain foods or drinks like fried foods and coffee.

This is often related to pregnancy sickness, but may also be the body's way of ensuring that they eat and drink wisely. Generally, there's no harm in giving into food cravings, especially if doing so helps getting through phases like early morning sickness which can be pretty distressing. However this must be done in moderation.

Eating a lot of one food only means eating less of other foods and therefore running the risk of becoming deficient in important nutrient. A craving for non-food items — such as ash or soap or toothpaste — is known as `pica'.Pica is potentially very harmful if indulged in and must be resisted.

Also, substances like soaps and ash may prevent the absorption of nutrients and other food substances. Courtesy : The Hindu

Why do doctors prescribe some medicines to be taken before and some after food?

When we take a medication, it is absorbed from various parts of our gut — some get absorbed in the stomach, some pass through the stomach into the intestines and get absorbed there.

The most important reason for timing a medication that is taken orally is to maximise its absorption so that more of the medication goes through the stomach into the blood.

Many medicines get absorbed better when food is not present and hence are taken on an empty stomach (an example is the hormone thyroxine which must be taken first thing in the morning). A few actually get absorbed better when food is present — an example of this is the antibiotic azithromycin. Some drugs are taken specifically with or after food because this may reduce the side effects of the drug on the stomach.

For example pain medications and certain antibiotics all can irritate the lining of the stomach and therefore are best taken with or after food. Some drugs work in the wall of the stomach to reduce the absorption of food and this is the desired therapeutic effect — an example of this is the anti-diabetic drug called acarbose.

This drug must be taken with the first bite of food. Similarly other oral anti-diabetic drugs and the injection insulin are taken before food because that is when they need to act — just after you eat a meal. The long and short of it is that the timing of a drug has important effects on its absorption, action, potency and even side effects and it is a good idea not to leave the doctor's office without being sure when you should be taking your medicine. Courtesy : The Hindu

A hen lays an egg every day. From where does it get the calcium required to make the eggshell?

 A commercial layer that lays an egg almost everyday derives its calcium requirements entirely from the feed it consumes. The nutritional requirements of chicken have been understood completely over the years. Today's layer bird yields 325 eggs in 365 days.

The chicken feed is mainly composed of maize, broken rice, soya, groundnut and sunflower cake, oyster shell grit and dicalcium phosphate.

The calcium requirements are mainly contributed by the shell grit. The feed is a precisely formulated one and the nutritional requirements of the layer are fully met by the feed it consumes daily. Courtesy : The Hindu

Why does it take a longer time to copy a file to a computer than delete it?

An Operating System (OS) is a program that acts as an intermediary between a user of a computer and the computer hardware.

Depending on the storage device being used, computers can store information in several physical forms. Each device has its own characteristics and physical organization, and hence different views of information are created. To unify all these views of information, a uniform logical view called a file is created.

A file is a contiguous set of data. It is the job of the OS to map this sequence of data into physical devices. The part of the OS responsible for this is the file system.

So the main task of the file system is to free the users of the details of storing of information in the physical devices. That is, when the storage device is changed, from disk to CD for example, the user still sees the same information.

In the most basic form, a file system consists of two distinct parts: a collection of files and a directory structure. The directory structure organizes and provides information about all the files in the system.

Every file has certain attributes like its name, location (its address in the file system), size, access control information (whether the file can be read or written to or only executed etc), type (whether it is just a collection of data or has other special instructions), time, date and user identification.

If file A is to be copied to file B, then first a new file called B is created, next the contents of A are read and finally, this is written to B.

For this to happen, three steps are necessary. First, space must be found for B in the file system. Second, an entry for B must be made in the directory. The directory records the name of B and the location of B in the file system. Third, a request is sent to the OS to read the contents of A.

The OS finds the location information of A from the directory and reads the contents. Now to write to B, the OS again searches the directory for the address of B. Finally the content of A is written onto the space provided for B.

On the other hand if a file deletion is to take place, the process is much simpler. If file A is to be deleted, the OS just searches the directory for the named file.

Having found the entry, the space occupied by A is released so that it can be used by other files and the directory entry is erased.

As can be seen, the operations involved while copying a file are much more than those involved while deleting a file. Courtesy : The Hindu

How are snake pits formed?

ANSWER I: Snake pits are a network of caves and crevasses formed by underground water and collapsed limestone serving as perfect locations for hibernating snakes. In cold countries, these pits protect them from very cold temperatures, which tend to dip as low as minus 40 degrees Celsius during winter. The snakes huddle themselves below the frost lines during harsh winter.

Snakes of the tropical countries, however protect themselves from the heat of summer by staying in pits, which are formed far below the ground level. These are protected underground from surface heat.

ANSWER II : Most snakes do not make their own burrows, but inhabit the burrows made by other creatures, such as rodents. They also inhabit termite mounds.

Only a few species of snakes, such as sandboas and shieldtail snakes, are capable of making their own burrows, but even they prefer existing burrows, when available - Courtesy : The Hindu

How are some insects able to walk on the surface of water?

 ANSWER I :Water has a see-through film on its top layer that is created by surface tension. That means molecules of water are more likely to cling to other molecules of water than to something else. Some insects have a waxy coating on their body/feet. The surface tension of this coating [20 to 30 ergs/cm{+2}] is much less than water [72 ergs/cm{+2}] and hence water tends to bond to itself rather than wetting the insect's feet and unless the insect is too heavy, it can remain on top of the water surface. For example, water striders and carpenter ants are so light that they can support themselves by spreading their weight on the surface tension of the water.

ANSWER II: Insects like the pond skater are able to walk on the surface of water mainly because of two factors. The first is their water-repellent (hydrophobic) cuticle. The cuticle of a pond skater is coated with wax to make it waterproof.

Again, the insects are able to maintain floatation or stand on the surface because a sufficiently large amount of its surface area is in contact with the water. The heavier the object, the more surface area is necessary to maintain floatation. As insects are very light weight the area of contact with the water surface is enough for it to prevent it from drowning. Courtesy : The Hindu

Why is it easier to tear wet paper and not dry paper?

Tearing a paper involves overcoming the cohesive force between the cellulose fibres (of which paper is made). In the case of dry paper this force is high and hence tearing it is not very easy. However, the cohesive force that is of electrostatic origin becomes weakened in the presence of water.

 

This is akin to the way table salt (sodium chloride) dissolves in water due the weakening of the electrostatic attraction between the positively and negatively charged ions. In the case of paper, the effect becomes easily perceptible as paper is hydrophilic and absorbs water.

Once dipped in water, the water molecules can easily flow into the spaces between the fibres, weakening the cohesive force between them and making them susceptible to easy tearing. Courtesy : The Hindu

How is the speed of a computer measured?

Two important factors that determine the speed of a computer are the amount of data that the Central Processing Unit can process in a given period of time and the CPU's clock speed.

The speed at which a CPU executes instructions is called the clock rate.

Every system contains an internal clock that regulates the rate at which instructions are executed and synchronizes all the various computer components. The CPU requires a fixed number of clock ticks to execute each instruction.

The faster the clock, the more instructions the CPU can execute per second. Clock speeds are expressed in megahertz MHz or gigahertz GHz. Mega means million and hertz means times per second, 200 MHz is 200 million times per second (and 200 GHz is 200 billion times per second).

The internal architecture of a CPU has as much to do with a CPU's performance as the clock speed. One common architecture is parallel processing. For example, while an instruction is being executed, the next instruction can be fetched from memory and decoded.

Instruction Prefetching is another idea where the CPU fetches the next instruction beforehand and places it in a queue for the execution unit to use the same.

The overall speed of a computer is also affected by the speed and size of the instruction/data bus. The instruction/data bus is the pathway for data communications between the computer's CPU and the various components in the computer.

The computer's bus has a certain size or width called the data path which is measured in bits and the speed of the bus is measured in MHz.

The larger the bus width and/or the faster the bus speed, the more data that can travel on it in a given amount of time.

Another factor affecting the speed is the size of the primary memory and cache. Increasing the size of the primary memory will speed up the performance if you run several applications at the same time or work with large files and documents. Cache is a small amount (normally less than 1 MegaByte) of high-speed memory residing on or close to the CPU. Cache memory supplies the CPU with the most frequently requested data and instructions.

Finally, effective interfacing of Input-Output devices to the CPU also increases the speed. Systems today use direct memory access (DMA) hardware wherein I/O device acts as a master and transfers large number of data to/from memory without intervention by the CPU. Courtesy : The Hindu

What is the powdery deposit found on some fruits like grapes? What is its use?

The white deposit seen on grapes and most other berries is cuticular wax. Cuticle is the outermost layer covering the plant surface and plays a role in the plant's interactions with its environment.Cuticular wax is part of the cuticle in several plant parts in almost all plant species. It is usually embedded in the cuticle and in some plant species crystalline wax structures overlay this layer and appear as powdery white/grey deposit (for example, grapes and other berries).

The wax is composed primarily of long-chain fatty acids, hydrocarbons, ketones, alcohols and alkaloids. Plants use cuticular wax primarily to regulate non-stomatal loss of water. 

Cuticular wax is also reported to play important roles in disease resistance against bacterial and fungal pathogens of plants and in plant-insect relationships.

Being reflective in nature, waxes are also thought to offer some protection against UV damage. In agriculture, waxes impede the uptake of foliar sprays without surfactants due their hydrophobicity (water repellent property). Courtesy : The Hindu

How do raw mangoes and bananas become ripe when treated with chemicals?

ANSWER I: The process of fruit ripening is chiefly regulated by a gaseous plant hormone called ethylene. Most fruits have elevated ethylene levels during ripening and sometimes just a peak in ethylene levels, just before the process of ripening begins.Ethylene regulates the expression of several genes involved in fruit ripening so as to modulate the activity of various enzymes involved in the process of ripening. These enzymes act to soften the `skin' of the fruit and also convert complex polysaccharides into simple sugars.

The chemical commonly used to ripen fruits commercially is ethephon (2-chloroethylphosphonic acid), which penetrates into the fruit and decomposes to ethylene. Incidentally, chemicals (e.g. calcium carbide) that produce acetylene, an analogue of ethylene, are also used in some places posing dangers of explosion and carryover of toxic materials to consumers.Ethylene is induced by several cues such as higher temperature, wounding, disease etc. Higher levels of ethylene and enhanced respiration might contribute to ripening when stored at higher temperatures.

ANSWER II: The ripening signal of a fruit comes form a hormone ethylene. Production of ethylene turns on some genes that are transcribed and translated to produce other enzymes. These enzymes are responsible for the conversion of starch into simple sugar, degradation of chlorophyll and appearance of other new pigments like carotenoids, change in the skin colour and the breakdown of acid, making the fruit taste neutral.

Hardy nature of the skin loosens when pectin is broken-down by an enzyme pectinase. Conversion of larger molecules into smaller volatile substances causes an aromatic odour.Natural process of fruit ripening is accelerated by using certain chemicals. Here, calcium carbide is used. When carbide is dissolved in water it produces acetylene, an analogue of ethylene, a natural fruit-ripening agent.

The ripening process is accelerated since acetylene imitates ethylene. Since the amount of carbide needed to ripen the immature fruit is more it makes the fruit become more tasteless and toxic. Presence of trace amount of arsenic and phosphorous in carbide makes the healthy fruits poisonous.One can distinguish the artificially ripenened fruit by the uniform skin colour in fruits like tomato, mango, papaws, etc and in the case of banana, yellow colour fruit with dark green stem. Courtesy : The Hindu

How does wheat flour become malleable and elastic when mixed with water?

 Basically wheat flour does not contain any malleable or elastic characteristic materials. When water is added to wheat flour, a new product called Gluten is formed by hydration of wheat proteins. It causes the production of dough. Gluten contains water approximately 2/3rd of its weight. It forms about 90 per cent of the total protein of flour. It is stretchable product just like rubber. It also contains small quantity of fat, cellulose and minerals.

 Gluten in turn contains protein fragments called Glutenin and Gliadin. These two confer the dough the elastic and malleable properties. When mixed together, as they are in dough, these two proteins form a tangle of strands that trap the gas. While gliadin in gluten confers mellowness and elasticity, the glutenin provides the structure. Greater the amount of gliadin, softer will be the gluten. Gluten is responsible for the rheological properties of dough because it forms the skeleton of the dough. High structured products like bread and bun require stronger quality of gluten while low structured products like biscuits and cakes do not require strong gluten. Courtesy : The Hindu

Why are we unable to see through a frosted glass?

A glass plate, which is polished on both the surfaces is `perfectly' transparent to light and so one can see through it. However if one of the surfaces is sand blasted to get a frost glass, this rough surface would scatter almost all the light in all directions.

Therefore light entering one side of the glass plate is totally scattered and lost and thus is not able to pass through. So we cannot see through. If one of the surfaces is mirrored then too one cannot see through but in this case the light is not scattered but completely reflected off and one would then see it as a shiny mirror. Courtesy : The Hindu

What causes milk to rise up when we boil it?

 Milk contains 87 per cent water, 4 per cent proteins and 5 per cent lactose (milk sugar). When we boil milk, the fat, sugar, proteins and minerals get separated. Since they are lighter than milk they collect on the surface in the form of cream.

During heating some amount of water gets converted into vapour and the bubbles of water vapour rise to the top but the heat is conducted away by the layer of water and by the fat droplets that have a higher boiling point than water.

The vapour gets trapped in the creamy upper layer. As the milk is heated further the water vapour expands and thick foam is produced on the top.

As the milk is boiled continuously the water, which boils at 100 degrees Centigrade, produces more water vapour and pressure builds up in the boiling milk so that the vapour pressure raises the creamy layer. So the milk pushes the creamy layer out and milk spills out. Courtesy : The Hindu

Why does an egg (with the shell) burst when cooked in a microwave oven?

 Microwave radiation is generated in an electronic tube called a magnetron, and passes along what's called a wave-guide into the oven cavity.The microwaves are absorbed by foods — a characteristic that make them ideal for cooking. The microwave energy transmitted in a microwave oven is directed toward the centre of the compartment. The highest absorption factor for microwave energy is water. The water absorbs the energy and becomes agitated and this molecular level agitation is the friction that heats up food

When microwaved, different components in an egg expand at different rates, which can result in the egg exploding. White portion of egg contains a high proportion of water and yolk contains a high proportion of fat. Microwaved eggs can reach temperatures much higher than if they were simply boiled in water at 100 degrees Celsius. At these elevated temperatures, water inside the egg, mostly in the white albumen, vapourises — even as the albumen solidifies. If the pressure inside the egg exceeds the breaking strength of the shell, the egg will explode.

Using a wooden pick or tip of a knife to break the yolk membrane of an unbeaten egg before micro cooking to allow the steam to escape, can help prevent the explosion. Covering cooking containers with a lid, plastic wrap or wax paper encourages even cooking and (if we forget to prick the yolk) helps to confine the explosion Courtesy : The Hindu

Why is fuel used in airplanes different from those used in motor vehicles?

 Aviation turbine fuels are used for powering jet and turbo-prop engine aircraft. Kerosene was used to fuel the first turbine engines. Kerosene-type fuel was chosen as having the best combination of properties. 

As the primary function of aviation turbine fuel (jet fuel) is to power an aircraft, energy content and combustion quality are key fuel performance properties. Other significant performance properties are stability, lubricity, fluidity, volatility, non-corrosivity, and cleanliness. Besides providing a source of energy, fuel is also used as a hydraulic fluid in engine control systems and as a coolant for certain fuel system components.

However, compared to a kerosene-type fuel, other type fuels like used in motor vehicles were found to have operational disadvantages due to their higher volatility:

• Greater losses due to evaporation at high altitudes
•  Greater risk of fire during handling on the ground
•  Crashes of planes fuelled with wide-cut fuel were less survivable
•  Lighter (less dense) fuels, such as gasoline, have higher heating values on a weight basis: whereas heavier (more dense) fuels, like diesel, have higher heating values on a volume basis. Since space is at a premium in most aircraft, the amount of energy contained in a give quantity of fuel is important. A fuel with high volumetric energy content maximises the energy that can be stored in a fixed volume and thus provides the longest flight range.

There are currently two main grades of turbine fuel in use in civil commercial aviation: jet A-1 and jet A, both are kerosene type fuels. There is another grade of jet fuel, jet B which is a wide cut kerosene (a blend of gasoline and kerosene) but it is rarely used except in very cold climates.

Jet A-1 is a kerosene grade of fuel suitable for most turbine engine aircraft. It is produced to a stringent internationally agreed standard, has a flash point above 38 degrees centigrade (100 degrees Fahrenheit) and a freeze point maximum of minus 47 degrees Centigrade.

Jet A is a similar kerosene type of fuel, produced and normally only available in the U.S. It has the same flash point as Jet A-1 but a higher maximum freeze point (minus 40 degrees centigrade).

Jet B is a distillate covering the naphtha and kerosene fractions. It can be used as an alternative to jet A-1 but because it is more difficult to handle (higher flammability), there is only significant demand in very cold climates where its better cold weather performance is important. Courtesy : The Hindu

How does a rechargeable battery work? What is the life of such batteries and how are they different from ordinary batteries?

 Electrochemical cells and batteries are identified generally as primary and secondary batteries. The primary batteries cannot be easily or effectively re-charged electrically and hence are discharged (used) and discarded. The electrochemical reactions in primary cells are not easily reversible. When the battery delivers current (during use) the active materials undergo changes and the active materials slowly will become inactive because the discharged active materials can't deliver current. In secondary batteries (example., lead-acid) the reactions are said to be reversible because once the battery is used, the inactive materials can be converted back to active materials by re-charging and the battery will be again ready for use.

These systems are also called as `storage batteries'. (example., lead-acid, nickel-cadmium) In the primary category, for example., zn-carbon cells, the anode is zinc and cathode is manganese dioxide. During discharge (when battery in use), the simplified reaction can be written as (the actual electrochemical process is more complicated)

Zn + 2 MnO{-2} ZnO + Mn{-2}O{-3}

Discharge (delivers current)

The discharged products (right hand side) cannot be formed back into original active materials (left hand side) by passing current in an opposite direction (charging). It is said to be `irreversible'

Where as in secondary batteries, for example., lead-acid, the active materials can be formed back after discharge (use) and it will be ready for use again after charge.

Pb + PbO{-2} + 2H{-2}SO{-4}

Technically some primary batteries can be recharged for several cycles but may not deliver full capacity and may have poor charge retention after recharge. Generally the cells are not designed for that type of use. The life of a secondary battery (lead-acid or nickel-cadmium) may vary from 200-1200 cycles (one cycle represents one discharge and charge) depending on its design parameters. Courtesy : The Hindu

Does frequent switching on/off of a fluorescent lamp reduce its life?

 The life of a fluorescent lamp is essentially determined by life of the cathode filament it uses. A conventional fluorescent lamp employs closely wound coil of tungsten wire as filament. Upon switching on the lamp, electric current passing through the filament will raise the temperature of the filament that in turn will generate thermions (electrons generated by a thermal process). Thermions are necessary to initiate electric-discharge through the column of the fluorescent lamp.

 Frequent switching on/off the fluorescent lamp occurs through several cycles of filament heating and cooling. If the cycles of heating and cooling of the filament are too frequent this may result in tremendous loss of oxide coating (at the rate of 10-20 micro-grams/cm{+2} per cycle).

  

 The loss of oxide coating in the cathode filament through rapid on/off (heating/cooling) operations will lead to poor performance of the filament in generating thermions to initiate the discharge process. This in turn will reduce the life of the fluorescent lamp. Life of a conventional fluorescent lamp usually rated for several thousand hours of continuous burning can be halved or made still less, just by frequent switching on/off. Courtesy : The Hindu

How do touch screens work?

 Touch screen monitors — where you can use your finger on the computer screen to navigate through the contents — have become more and more commonplace over the past decade, particularly at public information kiosks. A basic touch screen has three main components: a touch sensor, a controller, and a software driver. The touch screen is an input device, so it needs to be combined with a display and a PC to make a complete touch input system.

The Touch Sensor has a textured coating across the glass face. This coating is sensitive to pressure and registers the location of the user's finger when it touches the screen. The controller is a small PC card that connects the touch sensor to the PC. It takes information from the touch sensor and translates it into information that PC can understand. The Software Driver is a software update for the PC system that allows the touchscreen and computer to work together. It tells the computer's operating system how to interpret the touch event information that is sent from the controller.

There are three basic systems that are used to recognise a person's touch — Resistive, Capacitive and Surface acoustic wave.

The resistive system consists of a normal glass panel that is covered with a conductive and a resistive metallic layer. These layers are held apart by spacers, and a scratch-resistant layer is placed on top of the whole set up. An electrical current runs through the two layers while the monitor is operational. When a user touches the screen, the two layers make contact in that spot. The change in electrical field is noted and coordinates of the point of contact are calculated. Once the coordinates are known, a special driver translates the touch into something that the operating system can understand, much as a computer mouse driver translates a mouse's movements into a click or drag.

In the capacitive system, a layer that stores electrical charge is placed on the glass panel of the monitor. When a user touches the monitor with his or her finger, some of the charge is transferred to the user, so the charge on the capacitive layer decreases. This decrease is measured in circuits located at each corner of the monitor. The computer calculates, from the relative differences in charge at each corner, exactly where the touch event took place and then relays that information to the touch screen driver software. One advantage of the capacitive system is that it transmits almost 90 per cent of the light from the monitor, whereas the resistive system only transmits about 75 per cent. This gives the capacitive system a much clearer picture than the resistive system.

The surface acoustic wave system uses two transducers (one receiving and one sending) placed along the x and y axes of the monitor's glass plate. Also placed on the glass are reflectors — they reflect an electrical signal sent from one transducer to the other. The receiving transducer is able to tell if the wave has been disturbed by a touch event at any instant, and can locate it accordingly. The wave setup has no metallic layers on the screen, allowing for 100-percent light throughput and perfect image clarity. This makes the surface acoustic wave system best for displaying detailed graphics (both other systems have significant degradation in clarity).

Another area in which the systems differ is which stimuli will register as a touch event. A resistive system registers a touch as long as the two layers make contact, which means that it doesn't matter if you touch it with your finger or a rubber ball. A capacitive system, on the other hand, must have a conductive input, usually your finger, in order to register a touch. The surface acoustic wave system works much like the resistive system, allowing a touch with almost any object — except hard and small objects like a pen tip. (Source: www.howstuffworks.com and www.touchscreens.com ) Courtesy : The Hindu

Why is the Earth's core hot? What caused it to heat up? Is it still heating, or now cooling?

Scientists estimate that temperature at the Earth's core is about 5538{+0}C.

Much of the heat inside the Earth today comes from elements that were present when the planet was first formed billions of years ago. One theory is that radioactive decay of the primordial elements inside the Earth, U-238, Th-232, and U-235 and their radioactive products generate thermal energy (heat).

A nucleus — the central core of an atom — contains both protons and neutrons. Elements, such as the ones mentioned above, have a fixed number of protons but may exist with various numbers of neutrons. The sum of the protons and neutrons makes up the mass number of an element. Isotopes of an element have the same chemical properties but different weights (indicated by the mass number). Radioactive elements are isotopes with an unstable nucleus.

The isotopes decay by emitting energetic alpha and beta particles until stability is reached. Alpha particles are the nuclei of ordinary helium atoms, which consist of two protons and two neutrons. Beta particles are electrons or positrons. The half-life of an isotope is the amount of time it takes for half of the atoms to decay into a more stable form. Within the Earth, the released particles from the elements are slowed by friction through interaction with Earth material, thereby generating heat. The primordial radioactive elements have half-lives on the order of a billion years. Hence, since the Earth formed, their abundance is decreasing over time as a function of their half-life. Therefore, Earth's core is not heating up, it's cooling down. Courtesy : The Hindu

Why is the @ symbol used in an email address?

ANSWER I: The e-mail address generally has two parts, user id and the domain name. The @ (at) symbol is used to separate the user id from the domain name in the e-mail address. The e-mail address is in the form, userid@domainname.com (example : rajiv@hotmail.com).

The domain name is usually the name of the service provider and it cannot be changed. The user name can be changed by creating a new e-mail address. Some websites like www.myownemail.com allow the users to have a domain name of their choice like rajiv@rajiv.com or rajiv@britneyspears.com or rajiv@quackquack.com or any other name on earth. The significance of the @ symbol is that it separates the user id from the domain name.

ANSWER II: Email addresses are basically identifiers of users and are unique. They are usually composed of the following parts, namely: username@subdomain.domain. Example: rajesh_27@yahoo.com. These addresses can be compared to the addresses of the houses in a huge colony. Examples are Hotmail, Yahoo, Sify and AOL.The domains are in turn classified as sub-domains for the sake of clarity in large organisations. Some of the top-level domains (TLD's) are com, .edu, .net,.gov. The `@' symbol indicates that the user can be reached on the Internet by giving the email address(also called the User's Uniform Resource Locator (URL)). Courtesy : The Hindu